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Stable Diffusion XL

Run SDXL at the speed of light

Input

Fields marked * are required. Review the price before generating.

The prompt to use for generating the image. Be as descriptive as possible for best results.

The size of the generated image.

The number of images to generate.

Additional settings
Tips and requirements

The negative prompt to use. Use it to address details that you don't want in the image. This could be colors, objects, scenery and even the small details (e.g. moustache, blurry, low resolution).

The number of inference steps to perform.

The same seed and the same prompt given to the same version of Stable Diffusion will output the same image every time.

Tips and requirements

The CFG (Classifier Free Guidance) scale is a measure of how close you want the model to stick to your prompt when looking for a related image to show you.

The list of LoRA weights to use.

The list of embeddings to use.

If set to true, the prompt will be expanded with additional prompts.

The format of the generated image.

Price$0.60 / run

You sign up first. Your settings are kept.

API code for this input

Complete the required fields and fix input errors to see API code.

How to use Stable Diffusion XL

Getting started

Build illustration or photo-style concepts with a specific subject, setting and light.

Review your result

Check anatomy, repeated objects and whether a LoRA overwhelms the intended composition.

Inputs and settings

Start with prompt and one image. Compare guidance_scale and num_inference_steps independently; add loras only when you need their effect.

For source-based changes, choose image-to-image with image_url. Inpainting also requires mask_url; both expose strength. Output format is JPEG or PNG.

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